diff --git a/README.md b/README.md
index 8af4766f6744..5b3660349e90 100644
--- a/README.md
+++ b/README.md
@@ -76,15 +76,13 @@ You need to accept the model license before downloading or using the Stable Diff
 
 ```python
 # make sure you're logged in with `huggingface-cli login`
-from torch import autocast
 from diffusers import StableDiffusionPipeline
 
 pipe = StableDiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", use_auth_token=True)
 pipe = pipe.to("cuda")
 
 prompt = "a photo of an astronaut riding a horse on mars"
-with autocast("cuda"):
-    image = pipe(prompt).images[0]  
+image = pipe(prompt).images[0]  
 ```
 
 **Note**: If you don't want to use the token, you can also simply download the model weights
@@ -104,8 +102,7 @@ pipe = StableDiffusionPipeline.from_pretrained("./stable-diffusion-v1-4")
 pipe = pipe.to("cuda")
 
 prompt = "a photo of an astronaut riding a horse on mars"
-with autocast("cuda"):
-    image = pipe(prompt).images[0]  
+image = pipe(prompt).images[0]  
 ```
 
 If you are limited by GPU memory, you might want to consider using the model in `fp16` as 
@@ -123,8 +120,7 @@ pipe = pipe.to("cuda")
 
 prompt = "a photo of an astronaut riding a horse on mars"
 pipe.enable_attention_slicing()
-with autocast("cuda"):
-    image = pipe(prompt).images[0]  
+image = pipe(prompt).images[0]  
 ```
 
 Finally, if you wish to use a different scheduler, you can simply instantiate
@@ -149,8 +145,7 @@ pipe = StableDiffusionPipeline.from_pretrained(
 pipe = pipe.to("cuda")
 
 prompt = "a photo of an astronaut riding a horse on mars"
-with autocast("cuda"):
-    image = pipe(prompt).images[0]  
+image = pipe(prompt).images[0]  
     
 image.save("astronaut_rides_horse.png")
 ```
@@ -160,7 +155,6 @@ image.save("astronaut_rides_horse.png")
 The `StableDiffusionImg2ImgPipeline` lets you pass a text prompt and an initial image to condition the generation of new images.
 
 ```python
-from torch import autocast
 import requests
 import torch
 from PIL import Image
@@ -190,8 +184,7 @@ init_image = init_image.resize((768, 512))
 
 prompt = "A fantasy landscape, trending on artstation"
 
-with autocast("cuda"):
-    images = pipe(prompt=prompt, init_image=init_image, strength=0.75, guidance_scale=7.5).images
+images = pipe(prompt=prompt, init_image=init_image, strength=0.75, guidance_scale=7.5).images
 
 images[0].save("fantasy_landscape.png")
 ```
@@ -204,7 +197,6 @@ The `StableDiffusionInpaintPipeline` lets you edit specific parts of an image by
 ```python
 from io import BytesIO
 
-from torch import autocast
 import torch
 import requests
 import PIL
@@ -234,8 +226,7 @@ pipe = StableDiffusionInpaintPipeline.from_pretrained(
 pipe = pipe.to(device)
 
 prompt = "a cat sitting on a bench"
-with autocast("cuda"):
-    images = pipe(prompt=prompt, init_image=init_image, mask_image=mask_image, strength=0.75).images
+images = pipe(prompt=prompt, init_image=init_image, mask_image=mask_image, strength=0.75).images
 
 images[0].save("cat_on_bench.png")
 ```
@@ -258,7 +249,6 @@ If you want to run the code yourself 💻, you can try out:
 - [Text-to-Image Latent Diffusion](https://huggingface.co/CompVis/ldm-text2im-large-256)
 ```python
 # !pip install diffusers transformers
-from torch import autocast
 from diffusers import DiffusionPipeline
 
 device = "cuda"
@@ -270,8 +260,7 @@ ldm = ldm.to(device)
 
 # run pipeline in inference (sample random noise and denoise)
 prompt = "A painting of a squirrel eating a burger"
-with autocast(device):
-    image = ldm([prompt], num_inference_steps=50, eta=0.3, guidance_scale=6).images[0]
+image = ldm([prompt], num_inference_steps=50, eta=0.3, guidance_scale=6).images[0]
 
 # save image
 image.save("squirrel.png")
@@ -279,7 +268,6 @@ image.save("squirrel.png")
 - [Unconditional Diffusion with discrete scheduler](https://huggingface.co/google/ddpm-celebahq-256)
 ```python
 # !pip install diffusers
-from torch import autocast
 from diffusers import DDPMPipeline, DDIMPipeline, PNDMPipeline
 
 model_id = "google/ddpm-celebahq-256"
@@ -290,8 +278,7 @@ ddpm = DDPMPipeline.from_pretrained(model_id)  # you can replace DDPMPipeline wi
 ddpm.to(device)
 
 # run pipeline in inference (sample random noise and denoise)
-with autocast("cuda"):
-    image = ddpm().images[0]
+image = ddpm().images[0]
 
 # save image
 image.save("ddpm_generated_image.png")
diff --git a/src/diffusers/models/unet_2d_condition.py b/src/diffusers/models/unet_2d_condition.py
index dd9e2e570bf7..d510589404bc 100644
--- a/src/diffusers/models/unet_2d_condition.py
+++ b/src/diffusers/models/unet_2d_condition.py
@@ -266,7 +266,12 @@ def forward(
         timesteps = timesteps.expand(sample.shape[0])
 
         t_emb = self.time_proj(timesteps)
-        emb = self.time_embedding(t_emb.to(self.dtype))
+
+        # timesteps does not contain any weights and will always return f32 tensors
+        # but time_embedding might actually be running in fp16. so we need to cast here.
+        # there might be better ways to encapsulate this.
+        t_emb = t_emb.to(dtype=self.dtype)
+        emb = self.time_embedding(t_emb)
 
         # 2. pre-process
         sample = self.conv_in(sample)
diff --git a/src/diffusers/pipelines/README.md b/src/diffusers/pipelines/README.md
index 3462f5ff518d..71841e023372 100644
--- a/src/diffusers/pipelines/README.md
+++ b/src/diffusers/pipelines/README.md
@@ -86,15 +86,13 @@ logic including pre-processing, an unrolled diffusion loop, and post-processing
 
 ```python
 # make sure you're logged in with `huggingface-cli login`
-from torch import autocast
 from diffusers import StableDiffusionPipeline, LMSDiscreteScheduler
 
 pipe = StableDiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", use_auth_token=True)
 pipe = pipe.to("cuda")
 
 prompt = "a photo of an astronaut riding a horse on mars"
-with autocast("cuda"):
-    image = pipe(prompt).images[0]  
+image = pipe(prompt).images[0]  
     
 image.save("astronaut_rides_horse.png")
 ```
@@ -104,7 +102,6 @@ image.save("astronaut_rides_horse.png")
 The `StableDiffusionImg2ImgPipeline` lets you pass a text prompt and an initial image to condition the generation of new images.
 
 ```python
-from torch import autocast
 import requests
 from PIL import Image
 from io import BytesIO
@@ -129,8 +126,7 @@ init_image = init_image.resize((768, 512))
 
 prompt = "A fantasy landscape, trending on artstation"
 
-with autocast("cuda"):
-    images = pipe(prompt=prompt, init_image=init_image, strength=0.75, guidance_scale=7.5).images
+images = pipe(prompt=prompt, init_image=init_image, strength=0.75, guidance_scale=7.5).images
 
 images[0].save("fantasy_landscape.png")
 ```
@@ -148,7 +144,6 @@ The `StableDiffusionInpaintPipeline` lets you edit specific parts of an image by
 ```python
 from io import BytesIO
 
-from torch import autocast
 import requests
 import PIL
 
@@ -173,8 +168,7 @@ pipe = StableDiffusionInpaintPipeline.from_pretrained(
 ).to(device)
 
 prompt = "a cat sitting on a bench"
-with autocast("cuda"):
-    images = pipe(prompt=prompt, init_image=init_image, mask_image=mask_image, strength=0.75).images
+images = pipe(prompt=prompt, init_image=init_image, mask_image=mask_image, strength=0.75).images
 
 images[0].save("cat_on_bench.png")
 ```
diff --git a/src/diffusers/pipelines/stable_diffusion/README.md b/src/diffusers/pipelines/stable_diffusion/README.md
index 3a600c5859e9..45c4e4798cf7 100644
--- a/src/diffusers/pipelines/stable_diffusion/README.md
+++ b/src/diffusers/pipelines/stable_diffusion/README.md
@@ -59,15 +59,13 @@ pipe = StableDiffusionPipeline.from_pretrained("./stable-diffusion-v1-4")
 
 ```python
 # make sure you're logged in with `huggingface-cli login`
-from torch import autocast
 from diffusers import StableDiffusionPipeline
 
 pipe = StableDiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", use_auth_token=True)
 pipe = pipe.to("cuda")
 
 prompt = "a photo of an astronaut riding a horse on mars"
-with autocast("cuda"):
-    image = pipe(prompt).images[0]  
+image = pipe(prompt).sample[0]  
     
 image.save("astronaut_rides_horse.png")
 ```
@@ -76,7 +74,6 @@ image.save("astronaut_rides_horse.png")
 
 ```python
 # make sure you're logged in with `huggingface-cli login`
-from torch import autocast
 from diffusers import StableDiffusionPipeline, DDIMScheduler
 
 scheduler = DDIMScheduler(beta_start=0.00085, beta_end=0.012, beta_schedule="scaled_linear", clip_sample=False, set_alpha_to_one=False)
@@ -88,8 +85,7 @@ pipe = StableDiffusionPipeline.from_pretrained(
 ).to("cuda")
 
 prompt = "a photo of an astronaut riding a horse on mars"
-with autocast("cuda"):
-    image = pipe(prompt).images[0]  
+image = pipe(prompt).sample[0]  
     
 image.save("astronaut_rides_horse.png")
 ```
@@ -98,7 +94,6 @@ image.save("astronaut_rides_horse.png")
 
 ```python
 # make sure you're logged in with `huggingface-cli login`
-from torch import autocast
 from diffusers import StableDiffusionPipeline, LMSDiscreteScheduler
 
 lms = LMSDiscreteScheduler(
@@ -114,8 +109,7 @@ pipe = StableDiffusionPipeline.from_pretrained(
 ).to("cuda")
 
 prompt = "a photo of an astronaut riding a horse on mars"
-with autocast("cuda"):
-    image = pipe(prompt).images[0]  
+image = pipe(prompt).sample[0]  
     
 image.save("astronaut_rides_horse.png")
 ```
diff --git a/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion.py b/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion.py
index bc6ca1efbd9b..705b23d5fa7e 100644
--- a/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion.py
+++ b/src/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion.py
@@ -260,19 +260,20 @@ def __call__(
         # Unlike in other pipelines, latents need to be generated in the target device
         # for 1-to-1 results reproducibility with the CompVis implementation.
         # However this currently doesn't work in `mps`.
-        latents_device = "cpu" if self.device.type == "mps" else self.device
         latents_shape = (batch_size, self.unet.in_channels, height // 8, width // 8)
+        latents_dtype = text_embeddings.dtype
         if latents is None:
-            latents = torch.randn(
-                latents_shape,
-                generator=generator,
-                device=latents_device,
-                dtype=text_embeddings.dtype,
-            )
+            if self.device.type == "mps":
+                # randn does not exist on mps
+                latents = torch.randn(latents_shape, generator=generator, device="cpu", dtype=latents_dtype).to(
+                    self.device
+                )
+            else:
+                latents = torch.randn(latents_shape, generator=generator, device=self.device, dtype=latents_dtype)
         else:
             if latents.shape != latents_shape:
                 raise ValueError(f"Unexpected latents shape, got {latents.shape}, expected {latents_shape}")
-            latents = latents.to(latents_device)
+            latents = latents.to(self.device)
 
         # set timesteps
         self.scheduler.set_timesteps(num_inference_steps)
diff --git a/tests/test_pipelines.py b/tests/test_pipelines.py
index 78a22ec3138b..ca82cf5ce9b9 100644
--- a/tests/test_pipelines.py
+++ b/tests/test_pipelines.py
@@ -1214,6 +1214,37 @@ def test_stable_diffusion_memory_chunking(self):
         assert mem_bytes > 3.75 * 10**9
         assert np.abs(image_chunked.flatten() - image.flatten()).max() < 1e-3
 
+    @slow
+    @unittest.skipIf(torch_device == "cpu", "Stable diffusion is supposed to run on GPU")
+    def test_stable_diffusion_text2img_pipeline_fp16(self):
+        torch.cuda.reset_peak_memory_stats()
+        model_id = "CompVis/stable-diffusion-v1-4"
+        pipe = StableDiffusionPipeline.from_pretrained(
+            model_id, revision="fp16", torch_dtype=torch.float16, use_auth_token=True
+        ).to(torch_device)
+        pipe.set_progress_bar_config(disable=None)
+
+        prompt = "a photograph of an astronaut riding a horse"
+
+        generator = torch.Generator(device=torch_device).manual_seed(0)
+        output_chunked = pipe(
+            [prompt], generator=generator, guidance_scale=7.5, num_inference_steps=10, output_type="numpy"
+        )
+        image_chunked = output_chunked.images
+
+        generator = torch.Generator(device=torch_device).manual_seed(0)
+        with torch.autocast(torch_device):
+            output = pipe(
+                [prompt], generator=generator, guidance_scale=7.5, num_inference_steps=10, output_type="numpy"
+            )
+            image = output.images
+
+        # Make sure results are close enough
+        diff = np.abs(image_chunked.flatten() - image.flatten())
+        # They ARE different since ops are not run always at the same precision
+        # however, they should be extremely close.
+        assert diff.mean() < 2e-2
+
     @slow
     @unittest.skipIf(torch_device == "cpu", "Stable diffusion is supposed to run on GPU")
     def test_stable_diffusion_text2img_pipeline(self):